Probing Diffusion Denoising Dynamics for Contrastive Representation Learning

📅 2026-07-09
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🤖 AI Summary
This work addresses the challenge of efficiently adapting the denoising dynamics of a pre-trained diffusion model for discriminative representation learning while preserving its generative capabilities. The proposed approach treats noisy latent variables at different denoising timesteps as multi-view augmentations of the same image. By freezing the Stable Diffusion backbone and applying lightweight fine-tuning via LoRA, the method introduces, for the first time, a joint optimization framework combining noise-level contrastive learning with reconstruction loss. This enables synergistic learning between generative and discriminative objectives without requiring training from scratch. Experimental results demonstrate strong performance: a linear probe achieves 80.1% top-1 accuracy on ImageNet-1K, and unconditional image generation at 256×256 resolution yields an FID of 5.56.
📝 Abstract
Text-to-image diffusion models exhibit unprecedented generative capability and contain rich intermediate representations that can be useful for discriminative vision tasks. Motivated by this observation, we study a focused question: how can the denoising dynamics of a pretrained diffusion model be adapted to support discriminative representation learning while preserving its generative behavior under parameter-efficient updates? We present D$^3$CL as an investigation of this question. Our key observation is that noisy latents at different diffusion timesteps can be interpreted as stochastic views of the same underlying image, enabling a contrastive objective to be coupled with the standard denoising reconstruction loss. This formulation provides a simple way to probe the interaction between generative denoising and discriminative representation learning without training from scratch. To keep the adaptation lightweight, we apply LoRA updates to a pretrained Stable Diffusion backbone while freezing the original model parameters. D$^3$CL provides strong empirical evidence that reconstruction and noise-level contrastive objectives can be complementary: on ImageNet-1K, it obtains 80.1% linear-probing accuracy and an FID of 5.56 for $256 \times 256$ unconditional generation. Additional ablations on the design space suggest that the usefulness of diffusion features depends on where and how denoising states are sampled. These results establish D$^3$CL as a parameter-efficient adaptation framework for pretrained diffusion models, showing that noise-level contrastive learning can structure denoising representations for discriminative tasks while maintaining generative performance.
Problem

Research questions and friction points this paper is trying to address.

diffusion models
contrastive learning
representation learning
denoising dynamics
parameter-efficient adaptation
Innovation

Methods, ideas, or system contributions that make the work stand out.

diffusion models
contrastive learning
parameter-efficient adaptation
denoising dynamics
LoRA
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